Skip to the Main Content

Note:These pages make extensive use of the latest XHTML and CSS Standards. They ought to look great in any standards-compliant modern browser. Unfortunately, they will probably look horrible in older browsers, like Netscape 4.x and IE 4.x. Moreover, many posts use MathML, which is, currently only supported in Mozilla. My best suggestion (and you will thank me when surfing an ever-increasing number of sites on the web which have been crafted to use the new standards) is to upgrade to the latest version of your browser. If that's not possible, consider moving to the Standards-compliant and open-source Mozilla browser.

February 23, 2025

Potential Functions and the Magnitude of Functors 1

Posted by Tom Leinster

Next: Part 2

In the beginning, there were hardly any spaces whose magnitude we knew. Line segments were about the best we could do. Then Mark Meckes introduced the technique of potential functions for calculating magnitude, which was shown to be very powerful. For instance, Juan Antonio Barceló and Tony Carbery used it to compute the magnitude of odd-dimensional Euclidean balls, which turn out to be rational functions of the radius. Using potential functions allows you to tap into the vast repository of knowledge of PDEs.

In this post and the next, I’ll explain this technique from a categorical viewpoint, saying almost nothing about the analytic details. This is category theory as an organizational tool, used to help us understand how the various ideas fit together. Specifically, I’ll explain potential functions in terms of the magnitude of functors, which I wrote about here a few weeks ago.

Before I can describe this categorical viewpoint on potential functions, I have to explain what potential functions are in the magnitude context, and why they’re very useful. That’s what I’ll do today.

This part of the story is about metric spaces. For now I’ll assume they satisfy all the classical axioms, including symmetry of the metric, meaning that d(a,a)=d(a,a)d(a, a') = d(a', a) for all points aa and aa'. When metric spaces are viewed as enriched categories, symmetry isn’t automatic — but we’ll come to that next time.

A weighting on a finite metric space AA is a function w:Aw: A \to \mathbb{R} such that for all aAa \in A,

aAe d(a,a)w(a)=1. \sum_{a' \in A} e^{-d(a, a')} w(a') = 1.

Everyone who sees this formula for the first time asks where the exponential comes from. Ultimately it’s because of the enriched category viewpoint (which again we’ll come to next time), but the short story is that exponential is essentially the only reasonable function that converts addition into multiplication.

For simplicity, I’ll assume here that every finite metric space AA has a unique weighting, which I’ll call w Aw_A. Since the definition of weighting involves the same number of equations as unknowns, this is generically true (and it’s always true for subspaces of n\mathbb{R}^n), even though there are exceptions.

The magnitude of AA is

|A|= aAw A(a). |A| = \sum_{a \in A} w_A(a) \in \mathbb{R}.

That’s for finite metric spaces. To extend the definition to compact metric spaces, there are various ideas you might try. You could define the magnitude of a compact space as the supremum of all the magnitudes of its finite subspaces. Or, you could take an ever-denser sequence of finite subsets of your compact space, then define its magnitude to be the limit of the magnitudes of the approximating subsets. Or, you could try replacing the sums in the formulas above by integrals, somehow.

Mark Meckes showed that all these approaches are equivalent. They all give the same definition of the magnitude of a compact space. (At least, this is true subject to a condition called “positive definiteness” which I won’t discuss and which always holds for subspaces of n\mathbb{R}^n.)

How do we actually calculate magnitude, say for compact subspaces AA of n\mathbb{R}^n?

In principle, you already know how do it. You run through all the finite subsets of AA, calculate the magnitude of each using the definition above, and take the sup. The trouble is, this procedure is incredibly hard to work with. It’s completely impractical.

A slightly more practical approach is to look for a “weight measure”, that is, a signed measure ww on AA such that for all aAa \in A,

Ae d(a,a)dw(a)=1. \int_A e^{-d(a, a')}\, d w(a') = 1.

As Mark showed, if ww is a weight measure then the magnitude of AA is given by |A|=w(A)|A| = w(A). This is the analogue of the formula |A|= aw(a)|A| = \sum_a w(a) for finite spaces.

Example   Take an interval [c,d][c, d] in \mathbb{R}. It’s an elementary exercise to check that

12(δ c+λ [c,d]+δ d) \tfrac{1}{2}(\delta_c + \lambda_{[c, d]} + \delta_d)

is a weight measure on [c,d][c, d], where δ c\delta_c and δ d\delta_d denote the Dirac deltas at cc and dd, and λ [c,d]\lambda_{[c, d]} is Lebesgue measure on [c,d][c, d]. It follows that [c,d][c, d] is the total mass of this measure, which is

1+12(dc). 1 + \tfrac{1}{2}(d - c).

In other words, it’s 11 plus half the length of the interval.

The trouble with weight measures is that very few spaces have them. Even Euclidean balls don’t, beyond dimension one. It turns out that we need something more general than a measure, something more like a distribution.

In another paper, Mark worked out exactly what kind of measure-like or distribution-like thing a weighting should be. The answer is very nice, but I won’t explain it here, because I want to highlight the formal aspects above all.

So I’ll simply write the weighting on a metric space AA as w Aw_A, without worrying too much about what kind of thing w Aw_A actually is. All that matters is that it behaves something like a measure: we can pair it with “nice enough” functions ϕ:A\phi: A \to \mathbb{R} to get a real number, which I’ll write as

Aϕ(a)dw A(a) \int_A \phi(a) \, d w_A(a)

or simply

Aϕdw A. \int_A \phi \, d w_A.

And I’ll assume that all the spaces AA that we discuss do have a unique weighting w Aw_A.

That’s the background. But back to the question: how do we actually calculate magnitude?

First main idea   Don’t look at metric spaces AA in isolation. Instead, consider spaces embedded in some big space BB that we understand well.

Think of the big space BB as fixed. The typical choice is n\mathbb{R}^n. One sense in which we “understand” n\mathbb{R}^n well is that we have a weight measure on it: it’s just Lebesgue measure divided by a constant factor c nc_n. This follows easily from the fact that n\mathbb{R}^n is homogeneous. (It doesn’t matter for anything I’m going to say, but the constant c nc_n is ne xdx\int_{\mathbb{R}^n} e^{-\|x\|}\, d x, for which there’s a standard formula involving π\pis and factorials.)

The potential function of a subspace ABA \subseteq B is the function h A:Bh_A: B \to \mathbb{R} defined by

h A(b)= Ae d(a,b)dw A(a). h_A(b) = \int_A e^{-d(a, b)} \, d w_A(a).

By definition of weighting, h Ah_A has constant value 11 on AA. But outside AA, it could be anything, and it turns out that its behaviour outside AA is very informative.

Although w Aw_A is a measure-like thing on the subspace AA of BB, we typically extend it by 00 to all of BB, and then the definition of the potential function becomes

h A(b)= Be d(x,b)dw A(x) h_A(b) = \int_B e^{-d(x, b)} \, d w_A(x)

Examples   In all these examples, I’ll take the big space BB to be \mathbb{R}.

  • Let A={0}A = \{0\}. Then the weighting w Aw_A on AA is the Dirac delta at 00, and the potential function h A:h_A: \mathbb{R} \to \mathbb{R} is given by

    h(b)=e |b| h(b) = e^{-|b|}

    (bb \in \mathbb{R}), whose graph looks like this:

    graph of potential function of a singleton

  • Let A=[2,2]A = [-2, 2]. As we’ve already seen,

    w A=12(δ 2+λ [2,2]+δ 2), w_A = \tfrac{1}{2}(\delta_{-2} + \lambda_{[-2, 2]} + \delta_2),

    and a little elementary work then reveals that the potential function h A:h_A: \mathbb{R} \to \mathbb{R} is given by

    h A(b)={e (2b) if b2, 1 if 2b2, e (b2) if b2, h_A(b) = \begin{cases} e^{-(-2 - b)} &\text{if}   b \leq -2, \\ 1 &\text{if}   -2 \leq b \leq 2, \\ e^{-(b - 2)} &\text{if}   b \geq 2, \end{cases}

    which looks like this:

    graph of potential function of an interval

  • In the two examples we just did, the potential function of AA is just the negative exponential of the distance between AA and the argument bb. That’s not exactly coincidence: as we’ll see in the next part, the function just described corresponds to a strict colimit, whereas the potential function corresponds to a lax colimit. So it’s not surprising that they coincide in simple cases.

    But this only happens in the very simplest cases. It doesn’t happen for Euclidean balls above dimension 11, or even for two-point spaces. For example, taking the subset A={1,1}A = \{-1, 1\} of B=B = \mathbb{R}, we have

    h A(x)=e |x+1|+e |x1|1+e 2, h_A(x) = \frac{e^{-|x + 1|} + e^{-|x - 1|}}{1 + e^{-2}},

    which looks like this:

    graph of potential function of a two-point space

    whereas be d(A,b)b \mapsto e^{-d(A, b)} looks like this:

    graph of negative exponential distance of a two-point space

    Similar, but different!

And now we come to the:

Second main idea   Work with potential functions instead of weightings.

To explain what this means, I need to tell you three good things about potential functions.

  • The potential function determines the magnitude:

    |A|= Bh Adw B. |A| = \int_B h_A \, d w_B.

    At the formal level, proving this is a one-line calculation: substitute the definition of h Ah_A into the right-hand side and follow your nose.

    For example, we saw that n\mathbb{R}^n has weighting λ/c n\lambda/c_n, where λ\lambda is Lebesgue measure and c nc_n is a known constant. So

    |A|=1c n nh A(x)dx |A| = \frac{1}{c_n} \int_{\mathbb{R}^n} h_A(x) \, d x

    for compact A nA \subseteq \mathbb{R}^n. Here dxd x refers to ordinary Lebesgue integration.

    (For n\mathbb{R}^n, and in fact a bit more generally, this result appears as Theorem 4.16 here.)

  • You can recover the weighting from the potential function. So, you don’t lose any information by working with one rather than the other.

    How do you recover it? Maybe it’s easiest to explain in the case when the spaces are finite. If we write Z BZ_B for the B×BB \times B matrix (e d(b,b)) b,bB(e^{-d(b', b)})_{b', b \in B} then the definition of the potential function can be expressed very succinctly:

    h A=Z Bw A. h_A = Z_B w_A.

    Here h Ah_A and w Aw_A are viewed as column vectors with entries indexed by the points of BB. (For w Aw_A, those entries are 00 for points not in AA.) Assuming Z BZ_B is invertible, this means we recover w Aw_A from h Ah_A as Z B 1h AZ_B^{-1} h_A. And something similar is true in the non-finite setting.

    However, what really makes the technique of potential functions sing is that when B= nB = \mathbb{R}^n, there’s a much more explicit way to recover the weighting from the potential function:

    w A=(IΔ) (n+1)/2h A w_A = (I - \Delta)^{(n + 1)/2} h_A

    (up to a constant factor that I’ll ignore). Here II is the identity and Δ\Delta is the Laplace operator, i=1 n 2x i 2\sum_{i = 1}^n \frac{\partial^2}{\partial x_i^2}. This is Proposition 5.9 of Mark’s paper.

    How much more of this bullet point you’ll want to read depends on how interested you are in the analysis. The fundamental point is simply that (IΔ) (n+1)/2(I - \Delta)^{(n + 1)/2} is some kind of differential operator. But for those who want a bit more:

    • To make sense of everything, you need to interpret it all in a distributional sense. In particular, this allows one to make sense of the power (n+1)/2(n + 1)/2, which is not an integer if nn is even.

    • Maybe you wondered why someone might have proved a result on the magnitude of odd-dimensional Euclidean balls only, as I mentioned at the start of the post. What could cause the odd- and even-dimensional cases to become separated? It’s because whether or not (n+1)/2(n + 1)/2 is an integer depends on whether nn is odd or even. When nn is odd, it’s an integer, which makes (IΔ) (n+1)/2(I - \Delta)^{(n + 1)/2} a differential rather than pseudodifferential operator. Heiko Gimperlein, Magnus Goffeng and Nikoletta Louca later worked out lots about the even-dimensional case, but I won’t talk about that here.

    • Finally, where does the operator (IΔ) (n+1)/2(I - \Delta)^{(n + 1)/2} come from? Sadly, I don’t have an intuitive explanation. Ultimately it comes down to the fact that the Fourier transform of e e^{-\|\cdot\|} is ξ(1+ξ) (n+1)/2\xi \mapsto (1 + \|\xi\|)^{-(n + 1)/2} (up to a constant). But that itself is a calculation that’s really quite tricky (for me), and it’s hard to see anything beyond “it is what it is”.

  • The third good thing about potential functions is that they satisfy a differential equation, in the situation where our big space BB is n\mathbb{R}^n. Specifically:

    {h A1 on A (IΔ) (n+1)/2h A0 on  nA. \begin{cases} h_A \equiv 1 &\text{on}   A \\ (I - \Delta)^{-(n + 1)/2} h_A \equiv 0&\text{on}   \mathbb{R}^n \setminus A. \end{cases}

    Indeed, the definition of weighting implies that h A1h_A \equiv 1 on AA, and the “second good thing” together with the fact that w Aw_A is supported on AA give the second clause.

    Not only do we have a differential equation for h Ah_A, we also have boundary conditions. There are boundary conditions at the boundary of AA, because of something I’ve been entirely vague about: the functions we’re dealing with are meant to be suitably smooth. There are also boundary conditions at \infty, because our functions are also meant to decay suitably fast.

    Maybe Mark will read this and correct me if I’m wrong, but I believe there are exactly the right number of boundary conditions to guarantee that there’s (typically? always?) a unique solution. In any case, the following example — also taken from Mark’s paper — illustrates the situation.

Example   Let’s calculate the magnitude of a real interval A=[c,d]A = [c, d] using the potential function method.

Its potential function h Ah_A is a function \mathbb{R} \to \mathbb{R} such that h A1h_A \equiv 1 on [c,d][c, d] and h Ah A=0h_A - h''_A = 0 on the rest of the real line. That differential equation comes from taking (IΔ) (n+1)/2(I - \Delta)^{-(n + 1)/2} in the case n=1n = 1.

The functions ff satisfying f=ff'' = f are those of the form ke ±xk e^{\pm x} for some constant kk, and in our case we’re free to choose the constant and the ±\pm sign differently on the two connected components of A\mathbb{R} \setminus A. So there are lots of solutions. But h Ah_A is required to be continuous and to converge to 00 at ±\pm \infty, and that pins it down uniquely: the one and only solution is

h A(x)={e (cx) if xc, 1 if cxd, e (xd) if dx. h_A(x) = \begin{cases} e^{-(c - x)} &\text{if}   x \leq c, \\ 1 &\text{if}   c \leq x \leq d, \\ e^{-(x - d)} &\text{if}   d \leq x. \end{cases}

I showed you the graph of this potential function above, in the case where A=[2,2]A = [-2, 2].

So the magnitude of A=[c,d]A = [c, d] is

1c 1 h A(x)dx, \frac{1}{c_1} \int_\mathbb{R} h_A(x)\, d x,

where c 1c_1 is the constant e |x|dx=2\int_\mathbb{R} e^{-|x|} \, d x = 2. This gives the answer:

|A|=1+12(dc)=1+12length(A). |A| = 1 + \tfrac{1}{2}(d - c) = 1 + \tfrac{1}{2}length(A).

 

The crucial point is that in this example, we didn’t have to come up with a weighting on AA. The procedure was quite mechanical. And that’s the attraction of the method of potential functions.

Next time, I’ll put all this into a categorical context using the notion of the magnitude of a functor, which I introduced here.

Posted at February 23, 2025 8:21 AM UTC

TrackBack URL for this Entry:   https://golem.ph.utexas.edu/cgi-bin/MT-3.0/dxy-tb.fcgi/3590

13 Comments & 0 Trackbacks

Re: Potential Functions and the Magnitude of Functors 1

Nice post, and I really look forward to the next one!

Regarding that question about the uniqueness of the solution to the differential equation, Barceló and Carbery proved a uniqueness result as Proposition 2 in their paper, in odd dimensions and when AA is sufficiently nice. I don’t remember offhand whether more general versions are stated as explicitly in the (substantially more technical) papers by Gimperlein, Goffeng, and Louca.

Mainly, though, I want to give credit where it’s due and say that the idea of using the potential function (and even the notation hh that we’ve been using for it) are due to Charles Clum, who did a summer research project on magnitude with him in 2012 when he was an undergraduate at CWRU. (He later went on to get a PhD at Ohio State.) Charles introduced the potential function as a tool in trying to prove a conjecture he’d made about the magnitude of finite subsets of n\mathbb{R}^n. As it turned out, he instead found a counterexample to his conjecture, without using potential functions. But his original plan of proof seemed fruitful enough to me that I decided to try pushing the idea of potential functions further.

Posted by: Mark Meckes on February 23, 2025 3:16 PM | Permalink | Reply to this

Re: Potential Functions and the Magnitude of Functors 1

Thanks! I see that in the Barceló and Carbery paper now. And thanks too for the background.

Posted by: Tom Leinster on February 23, 2025 8:43 PM | Permalink | Reply to this

Re: Potential Functions and the Magnitude of Functors 1

There’s a simple “soft” proof of the fact that e ξe^{-\Vert \xi \Vert} has Fourier transform c n(1+x 2) (n+1)/2c_n (1 + \Vert x \Vert^2)^{-(n+1)/2} (the latter is called the multivariate Cauchy distribution, by the way).

Let’s define the Cauchy distribution μ t\mu_t with parameter tt as the distribution of the random point where the (n+1)(n+1)-dimensional Brownian motion starting at 00 hits a hyperplane at distance tt from the origin. Observe the following facts:

  1. By the scale invariance of the Brownian motion, μ t\mu_t is just μ 1\mu_1 scaled by tt.

  2. Furthermore, the Brownian motion is conformally invariant up to a time change. Since the hyperplane can be mapped conformally to a sphere, with the starting point at its center, the Cauchy distribution is in fact the conformal image of the uniform measure on the sphere. Hence the explicit formula for its density.

  3. Cauchy distributions form a convolution semigroup, i.e. μ t+s=μ t*μ s\mu_{t+s} = \mu_t \ast \mu_s. This is because to hit a plane at distance t+st+s the Brownian motion has to hit a parallel plane at distance tt first, and then starting from there, hit the plane at distance ss (formally, this is a strong Markov property argument).

11 and 33 combined imply that its Fourier transform satisfies μ^(tξ)=(μ^(ξ)) t\hat \mu(t \xi) = (\hat \mu(\xi))^t. Together with the orthogonal invariance this means that μ^(ξ)=e constξ\hat \mu(\xi) = e^{- \mathrm{const} \cdot \Vert \xi \Vert}.

Posted by: Alexander Shamov on February 24, 2025 7:39 AM | Permalink | Reply to this

Re: Potential Functions and the Magnitude of Functors 1

And to fix the constant we can relate the exponential rate of decay of μ^\hat \mu to the distance to the closest singularity of the analytic continuation of μ\mu.

Posted by: Alexander Shamov on February 24, 2025 7:46 AM | Permalink | Reply to this

Re: Potential Functions and the Magnitude of Functors 1

That’s great! Thanks for such a clear and gentle explanation.

I was taught that the (1-dimensional) Cauchy distribution is what you get when you stand a certain distance from an infinite wall and throw a dart at it at an angle chosen uniformly at random (within the 180 degree range where it will actually hit the wall). The landing points are Cauchy distributed along the wall.

But you’re saying, among other things, that you also get a Cauchy distribution if instead of the darts flying in a straight line, they wander about in a Brownian way. I guess this must be related to the property you state in (3), since we can approximate Brownian motion by lots of small straight line segments joined together. Is that right?

By the way, here’s the tricky proof I had in mind, from Stein and Weiss’s book Introduction to Fourier Analysis on Euclidean Spaces.

Posted by: Tom Leinster on February 24, 2025 5:04 PM | Permalink | Reply to this

Re: Potential Functions and the Magnitude of Functors 1

I’m not sure I’m following the argument with small increments. To approximate the Brownian motion you need to keep them small, and definitely not extend them until they hit the wall.

But there’s something else. Throwing the dart at a uniformly random direction means mapping the uniform measure on the circle to the line by the “central projection”. This is not quite the same as the stereographic projection, which would be the conformal map I was talking about. But in dimension 1 we get the same Cauchy distribution from both projections! This is related to the well-known “double angle” theorem from high school geometry. But the corresponding statement for “angular areas” fails in higher dimensions, so in general the central and stereographic projections give different distributions.

Posted by: Alexander Shamov on February 25, 2025 7:02 AM | Permalink | Reply to this

Re: Potential Functions and the Magnitude of Functors 1

I’m not sure I’m following the argument with small increments.

I didn’t have an argument: I was just guessing how it might work, apparently wrongly :-)

In any case, I’d like to understand why the following two probability two distributions on \mathbb{R} are the same:

  1. Stand a certain distance from a wall and throw an object (which moves in a straight line) in a direction chosen uniformly at random.

  2. Stand a certain distance from a wall and release an object, which wanders about in Brownian motion.

The wall stretches off to infinity in both directions. The distribution on \mathbb{R} is given by where the object hits the wall. In both cases, it’s conditional on the object actually hitting the wall (which you didn’t mention explicitly in your original comment, but I guess is implicitly assumed — unless it almost surely happens anyway?).

If I understand right, these two distributions on \mathbb{R} are the same, with the proviso that “a certain distance” might not mean the same thing in 1 and 2. Description 1 is what I think of as the standard description of the Cauchy distribution. Description 2 is the one in your comment (for n=1n = 1).

Is there an easy way to see that these two distributions are the same?

And are they also the same in higher dimensions?

Posted by: Tom Leinster on February 25, 2025 8:00 PM | Permalink | Reply to this

Re: Potential Functions and the Magnitude of Functors 1

Well, as I said in the previous comment, they’re not the same, except in dimension 1. And what happens in dimension 1 is that choosing the direction uniformly at random is the same thing as going towards a uniformly-random point on a circle that contains the starting point (as opposed to being centered at it).

As for the object hitting the wall - yes, it happens almost surely.

Posted by: Alexander Shamov on February 27, 2025 10:27 AM | Permalink | Reply to this

Re: Potential Functions and the Magnitude of Functors 1

Thanks again. As you will have detected, there’s quite a bit here that I don’t understand, but I do appreciate your explanations.

I think my take-home lesson is that although the description of the 1-dimensional Cauchy distribution in terms of straight-flying darts is correct, it’s in a sense misleading, as in higher dimensions the more interesting trajectory to consider is Brownian motion.

Zooming out, the whole point of this was to try to get some intuition for why the (pseudo)differential operator (IΔ) (n+1)/2(I - \Delta)^{(n + 1)/2} appears in the theory of magnitude. I mentioned in the post that this is ultimately because (1+ξ 2) (n+1)/2(1 + \|\xi\|^2)^{-(n + 1)/2} is the Fourier transform of x 2\|x\|_2. You’ve given a nice explanation for this fact about Fourier transforms, so I guess the challenge now is to use that to build intuition about this differential operator.

Posted by: Tom Leinster on February 27, 2025 3:44 PM | Permalink | Reply to this

Re: Potential Functions and the Magnitude of Functors 1

Not related to the specific content of this post, but seeing those e d(a,a)e^{-d(a, a')} put me in mind of my sole contribution to physics, which has us show a certain distance kernel to be positive semi-definite.

This is defined on symmetric groups where the distance between two elements is a fixed multiple of the minimal number of transpositions between them. Then e βd(σ,τ)e^{-\beta \cdot d(\sigma, \tau)} is positive semi-definite when e βe^{\beta} is a natural number.

Has anyone looked at the magnitude of metric spaces on groups?

Posted by: David Corfield on February 25, 2025 5:13 PM | Permalink | Reply to this

Re: Potential Functions and the Magnitude of Functors 1

A little bit. This is something that comes up every now and again: stuff that geometric group theorists do with word metrics seems like it has something in common with magnitude. It’s the kind of question that gets asked after colloquia, and it does seem promising, but I’m not aware that anyone’s really got anywhere with it.

Posted by: Tom Leinster on February 25, 2025 8:03 PM | Permalink | Reply to this

Re: Potential Functions and the Magnitude of Functors 1

For anyone who wants to do this computationally using a Cayley graph and knows MATLAB, the code at https://mathoverflow.net/a/391977/1847 could be useful.

Posted by: Steve Huntsman on February 26, 2025 1:17 PM | Permalink | Reply to this

Re: Potential Functions and the Magnitude of Functors 1

I said almost nothing in my post about why it’s the case that the potential function hh of a subset AA of n\mathbb{R}^n satisfies the (pseudo)differential equation

(IΔ) (n+1)/2h=0 (I - \Delta)^{(n + 1)/2} h = 0

on nA\mathbb{R}^n \setminus A. I skipped this explanation because I don’t have anything new or categorical to say about it. My only little hint was that it had to do with the Fourier transform of xe xx \mapsto e^{-\|x\|} being ξ(1+ξ 2) (n+1)/2\xi \mapsto (1 + \|\xi\|^2)^{-(n + 1)/2}. But here seems as good a place as any to tell a bit more of the story, which I feel like doing after Alexander was so helpful in explaining why that Fourier transform is what it is.

In keeping with the spirit of these posts, I’ll suppress all the serious analysis and just push symbols around, not worrying about what anything means or whether anything’s well-defined. I’ll also skip constant real factors in various places.

We start with a subset AA of n\mathbb{R}^n. The weighting on AA is written as w Aw_A, and can be seen as a measure (or measure-like thing) on n\mathbb{R}^n by extending by 00. The potential function h A:Bh_A: B \to \mathbb{R} is defined by

h A(b)= ne abdw A(b). h_A(b) = \int_{\mathbb{R}^n} e^{-\|a - b\|} \, d w_A(b).

If we write E: nE: \mathbb{R}^n \to \mathbb{R} for the function xe xx \mapsto e^{-\|x\|} then this can be written snappily as the convolution

h A=E*w A. h_A = E * w_A.

We now want to rearrange this equation to make w Aw_A the subject. In other words, we want an inverse to the process of convolving with EE. And Fourier transforms are how we get it.

Fourier transforms turn convolutions into products, so

h A^=E^w A^. \widehat{h_A} = \widehat{E} \cdot \widehat{w_A}.

And we know what E^\widehat{E} is: it’s ξ(1+ξ 2) (n+1)/2\xi \mapsto (1 + \|\xi\|^2)^{-(n + 1)/2}. So

w A^(ξ)=(1+ξ 2) (n+1)/2h A^(ξ). \widehat{w_A}(\xi) = (1 + \|\xi\|^2)^{(n + 1)/2} \cdot \widehat{h_A}(\xi).

One of the first things one learns about one-dimensional Fourier transforms is that the Fourier transform of a derivative ff', evaluated at ξ\xi, is iξf^(ξ)i\xi \hat{f}(\xi). (Hence the usefulness of Fourier transforms in solving differential equations.) Taking a couple of further steps down this path, we get a formula for the Fourier transform of the Laplacian Δf\Delta f of a function ff on n\mathbb{R}^n: it’s ξ 2f^(ξ)-\|\xi\|^2 \hat{f}(\xi). Since the Fourier transform is linear, it follows that

(1+ξ 2) (n+1)/2h A^=((IΔ 2) (n+1)/2h A) (ξ), (1 + \|\xi\|^2)^{(n + 1)/2} \cdot \widehat{h_A} = ((I - \Delta^2)^{(n + 1)/2} h_A)^{\wedge}(\xi),

at least when nn is odd. (The point is that when nn is odd, (1+ξ 2) (n+1)/2(1 + \|\xi\|^2)^{(n + 1)/2} is a polynomial in ξ 2\|\xi\|^2.) Hence

w A^=((IΔ 2) (n+1)/2h A) . \widehat{w_A} = ((I - \Delta^2)^{(n + 1)/2} h_A)^{\wedge}.

Finally, taking inverse Fourier transforms on each side gives

w A=(IΔ 2) (n+1)/2h A. w_A = (I - \Delta^2)^{(n + 1)/2} h_A.

And since w Aw_A is supported on AA, it follows that (IΔ 2) (n+1)/2h A0(I - \Delta^2)^{(n + 1)/2} h_A \equiv 0 on nA\mathbb{R}^n\setminus A.

Although this argument might look highly specific to n\mathbb{R}^n and the base of enrichment V=[0,]V = [0, \infty], and some of it genuinely is specific to this case, the overall shape of the argument is somewhat familiar if you’ve encountered Möbius inversion of some kind. Taking the equation

h A=E*w A h_A = E * w_A

and rearranging to get w Aw_A in terms of h Ah_A is a version of Möbius inversion. Another way to write that equation is h A=Zw Ah_A = Z w_A, where ZZ is the kernel (a,b)e ab(a, b) \mapsto e^{-\|a - b\|} and the product on the right-hand side is something like matrix multiplication.

The inverse operator Z 1Z^{-1} would be called the Möbius function in another context, and the observation that it’s something like a differential operator is nothing new. Indeed, when Rota wrote his first paper on Möbius inversion on posets, he described it as a “difference calculus” and wrote:

The [Möbius] inversion can be carried out by the analog of the “difference operator” relative to a partial ordering.

(Gian-Carlo Rota, On the foundations of combinatorial theory I: theory of Möbius functions, p.341).

Posted by: Tom Leinster on March 1, 2025 10:32 PM | Permalink | Reply to this

Post a New Comment